An update to my previous post on Stable Cascade with ComfyUI - instead of requiring four separate model files, we now only need two checkpoints, and the ComfyUI workflow is now very straightfoward!

On 12 Feb 2024, Stability.ai released Stable Cascade “research preview” (non-commercial license), and over the weekend, ComfyUI was updated to support this new model! Time to give it a go!

As a follow up to my last post regarding Consistent portraits using IP-Adapters for SDXL, this is a short comparison of the two face IP-Adapters for SDXL by h94 / xiaohu: namely, ip-adapter-plus-face_sdxl_vit-h.bin and ip-adapter-faceid_sdxl.bin.

Getting consistent character portraits generated by SDXL has been a challenge... until now! ComfyUI IPAdapter Plus (dated 30 Dec 2023) now supports both IP-Adapter and IP-Adapter-FaceID (released 4 Jan 2024)!

In the span of a couple of weeks, we got Crazy fast image generation with LCM LoRA for SDXL, which led me to ask if I could get Faster Stable Diffusion on M-series macs?. A few hours ago, Stability.ai gave us their response in the form of SDXL-Turbo... and now we go even faster!

All my recent Stable Diffusion XL experiments have been on my Windows PC instead of my M2 mac, because it has a faster Nvidia 2060 GPU with more memory. But today, I’m curious to see how much faster diffusion has gotten on a M-series mac (M2 specifically).

Time to try another ControlNet for Stable Diffusion XL - QR Code Monster v1 in ComfyUI. This ControlNet can influence SDXL such that the generated image “hides” a scan-able QR code, which at first glance, looks like a photo!

I never tried generating video clips or animations with SDXL before, simply because my GPU is not powerful enough. But after testing out the LCM LoRA for SDXL yesterday, I thought I’d try the SDXL LCM LoRA with Hotshot-XL, which is something akin to AnimateDiff.

Stable Diffusion keeps improving at an astounding pace! This time, it’s the idea of distilling a model into a Latent Consistency Model (LCM) for very, very fast image generation with a quality trade-off. On 24 Oct 2023, the distilled Segmind Stable Diffusion 1B (SSD-1B) model was released, followed by a better implementation in the form of Latent Consistency LoRAs for SDXL and SDD-1B released on 9 Nov 2023.

Stability AI just released an new SD-XL Inpainting 0.1 model. Here is how to use it with ComfyUI.

Stable Diffusion XL has trouble producing accurately proportioned faces when they are too small. Today, I learn to use the FaceDetailer and Detailer (SEGS) nodes in the ComfyUI-Impact-Pack to fix small, ugly faces.

I previously tried Thibauld’s SDXL-controlnet: OpenPose (v2) ControlNet in ComfyUI with poses either downloaded from OpenPoses.com or created with OpenPose Editor. Here are a few more options for anyone looking to create custom poses.

Yesterday, I tried out Stability AI’s four Control-LoRAs but mentioned that I did not understand the output of the Revision “image-mixing” workflow. I’ve since done a bit more experimentation...

Less than a week after my post testing diffusers/controlnet-canny-sdxl-1.0, along comes Stability AI’s own ControlNets, which they call Control-LoRAs! Not one but 4 of them - Canny, Depth, Recolor and Sketch models!

As promised in my last post, today I am testing out Thibaud Zamora’s SDXL-controlnet: OpenPose (v2) model using ComfyUI. I keep saying I’ll keep my posts short but never do...

Recently I tried Fooocus by Lyumin Zhang (Illyasviel) which fulfills its promise to allow one to “Focus on prompting and generating” - it is certainly easy to use! But shortly after its release, someone has “ported” the code to ComfyUI as a Custom Node! So of course it’s time to test it out...

The Stability AI documentation now has a pipeline supporting ControlNets with Stable Diffusion XL! Time to try it out with ComfyUI for Windows.

In this post, I experiment with latent scaling and latent compositing with SDXL 1.0 using ComfyUI. That is to say, increasing / decreasing the size of the image, and combining multiple images into one à la green screen (chroma key) compositing.

SDXL 1.0 uses two text prompts used to guide image generation. In my first post, SDXL 1.0 with ComfyUI, I referred to the second text prompt as a “style” but I wonder if I am correct. I have no idea! So let’s test out both prompts...

Onward with SDXL and ComfyUI! Sometimes I want to tweak generated images by replacing selected parts that don’t look good while retaining the rest of the image that does look good. Rather than manually creating a mask, I’d like to leverage CLIPSeg to generate a masks from a text prompt.