I’ve been playing around with the Stable Diffusion scripts a little (to be exact, Ben Firshman’s version). To help me understand the script, I decided to re-write it the way I prefer to use it... either breaking or optimizing it in the process :P

This is post number three for Stable Diffusion! Previously: 1. AI-generated images with Stable Diffusion on an M1 mac and 2. Stable Diffusion image-to-image mode.

My simplified code

I merged and simplified the original txt2img.py and img2img.py scripts... though I still don’t understand anything!

I removed:

  • input arguments - I find it easier just to edit the script and then run it,
  • prompt file,
  • grid output,
  • check_safety() - though the Stable Diffusion pipeline still will blank out NSFW images (don’t be surprised if you get black image),
  • put_watermark() - I can’t seem to read back the watermark anyway,
  • fallback model,
  • tqdm progress bars, and
  • transformer warning messages.

I hardcoded:

  • config YAML file,
  • checkpoint file,
  • use of PLMSSampler with ddim_eta=0 for txt2img, and use of DDIMSampler for img2img,
  • latent channels - which cannot be changed anyway, and
  • precision_scope based on device - MPS is f32 only anyway.

I changed:

  • seed_everything() to occur after the model is loaded when FIXED=1; otherwise, like in the original code, seed_everything() happens before the model is loaded (admittedly I’m uncertain about start_code), and
  • the output filename, so as to avoid filename collisions.

I added:

  • timers,
  • negative prompts - which seem to work great for me, and
  • a history file so that I can recall the prompt and seed.

If you don’t know what you are doing, please don’t just run my code. I consider this blog my diary... I do not release production quality code, and I barely test or handle errors in my code, dear reader.

The pre-requisite is a working, local copy of Stable Diffusion. For my setup, see my previous post, AI-generated images with Stable Diffusion on an M1 mac and the followup, Stable Diffusion image-to-image mode.


  • Place this script in the base stable-diffusion folder (not in the scripts folder),
  • Make sure to switch to the Python virtual environment e.g. source venv/bin/activate,
  • Mark the script as executable, chmod +x sd.py,
  • Run the script directly from the folder, ./sd.py. Note: My hashbang #! does not specify the exact python executable - this is not good practice! However, on my mac, since a python command is not found unless I ran setup the virtual environment above, this is good enough for me.
#! python
# copyright (c) 2022 C.Y. Wong, myByways.com simplified Stable Diffusion v0.1

import os, sys, time
import torch
import numpy as np
from omegaconf import OmegaConf
from PIL import Image
from einops import rearrange
from pytorch_lightning import seed_everything
from contextlib import nullcontext
from ldm.util import instantiate_from_config
from ldm.models.diffusion.plms import PLMSSampler
from ldm.models.diffusion.ddim import DDIMSampler
from transformers import logging

PROMPTS = [         # --prompt, one or more in an array
    'A high definition cartoon of humans fighting aliens on mars',
NEGATIVES = [       # negative prompt, one or more, default None (or an empty array)

HEIGHT = 512        # --H, default 512, beyond causes M1 to crawl
WIDTH = 512         # --W, default 512, beyond causes M1 to crawl
FACTOR = 8          # --f downsampling factor, default 8

FIXED = 0           # --fixed_code, 1 for repeatable results, default 0
SEED = 42           # --seed, default 42
NOISE = 0.0         # --ddim_eta, 0 deterministic, no noise - 1.0 random noise, ignored for PLMS (must be 0)
PLMS = 0            # --plms, default 1 on M1 for txt2img but ignored for img2img (must be DDIM)
ITERATIONS = 1      # --n_iter, default 1
SCALE = 7.5         # --scale, 5 further -> 15 closer to prompt, default 7.5
STEPS = 50          # --ddim_steps, practically little improvement >50 but takes longer, default 50

IMAGE = None        # --init-img, img2img initial latent seed, default None
STRENGTH = 0.75     # --strength 0 more -> 1 less like image, default 0.75

FOLDER = 'outputs'  # --outdir for images and history file below
HISTORY = 'history.txt'
CONFIG = 'configs/stable-diffusion/v1-inference.yaml'
CHECKPOINT = 'models/ldm/stable-diffusion-v1/model.ckpt'

def seed_pre():
    if not FIXED:

def seed_post(device):
    if FIXED:
        return torch.randn([1, 4, HEIGHT // FACTOR, WIDTH // FACTOR], device='cpu').to(torch.device(device.type))
    return None

def load_model(config, ckpt=CHECKPOINT):
    pl_sd = torch.load(ckpt, map_location='cpu')
    sd = pl_sd['state_dict']
    model = instantiate_from_config(config.model)
    model.load_state_dict(sd, strict=False)
    return model

def set_device(model):
    if torch.backends.mps.is_available():
        device = torch.device('mps')
        precision = nullcontext
    elif torch.cuda.is_available():
        device = torch.device('cuda')
        precision = torch.autocast
        device = torch.device('cpu')
        precision = torch.autocast
    return device, precision

def load_image(image_file):
    image = Image.open(image_file).convert('RGB')
    w, h = image.size
    w, h = map(lambda x: x - x % 32, (w, h))
    image = image.resize((w, h), resample=Image.Resampling.LANCZOS)
    image = np.array(image).astype(np.float32) / 255.0
    image = image[None].transpose(0, 3, 1, 2)
    image = torch.from_numpy(image)
    return 2.*image - 1.0

def setup_sampler(model):
    global NOISE
    if IMAGE:
        image = load_image(IMAGE).to(model.device.type)
        init_latent = model.get_first_stage_encoding(model.encode_first_stage(image))
        sampler = DDIMSampler(model)
        sampler.make_schedule(ddim_num_steps=STEPS, ddim_eta=NOISE, verbose=False)
        t_enc = int(STRENGTH * STEPS)
        sampler.t_enc = t_enc
        sampler.z_enc = sampler.stochastic_encode(init_latent, torch.tensor([t_enc]).to(model.device.type))
    elif PLMS:
        sampler = PLMSSampler(model)
        NOISE = 0
        sampler = DDIMSampler(model)
    return sampler

def get_prompts():
    global NEGATIVES
    if NEGATIVES is None:
        NEGATIVES = [''] * len(PROMPTS)
        NEGATIVES.extend([''] * (len(PROMPTS)-len(NEGATIVES)))
    return zip(PROMPTS, NEGATIVES)

def generate_samples(model, sampler, prompt, negative, start):
    uncond = model.get_learned_conditioning(negative) if SCALE != 1.0 else None
    cond = model.get_learned_conditioning(prompt)
    if IMAGE:
        samples = sampler.decode(sampler.z_enc, cond, sampler.t_enc, 
            unconditional_guidance_scale=SCALE, unconditional_conditioning=uncond)
        shape = [4, HEIGHT // FACTOR, WIDTH // FACTOR]
        samples, _ = sampler.sample(S=STEPS, conditioning=cond, batch_size=1,
            shape=shape, verbose=False, unconditional_guidance_scale=SCALE, 
            unconditional_conditioning=uncond, eta=NOISE, x_T=start)
    x_samples = model.decode_first_stage(samples)
    x_samples = torch.clamp((x_samples + 1.0) / 2.0, min=0.0, max=1.0)
    return x_samples

def save_image(image):
    name = f'{time.strftime("%Y%m%d_%H%M%S")}.png'
    image = 255. * rearrange(image.cpu().numpy(), 'c h w -> h w c')
    img = Image.fromarray(image.astype(np.uint8))
    img.save(os.path.join(FOLDER, name))
    return name

def save_history(name, prompt, negative):
    with open(os.path.join(FOLDER, HISTORY), 'a') as history:
        history.write(f'{name} -> {"PLMS" if PLMS else "DDIM"}, Seed={SEED}{" fixed" if FIXED else ""}, Scale={SCALE}, Steps={STEPS}, Noise={NOISE}')
        if IMAGE:
            history.write(f', Image={IMAGE}, Strength={STRENGTH}')
        if len(negative):
            history.write(f'\n + {prompt}\n - {negative}\n')
            history.write(f'\n + {prompt}\n')

def main():
    print('*** Loading Stable Diffusion - myByways.com simple-sd version 0.1')
    tic1 = time.time()
    os.makedirs(FOLDER, exist_ok=True)

    config = OmegaConf.load(CONFIG)
    model = load_model(config)
    device, precision_scope = set_device(model)
    sampler = setup_sampler(model)
    start_code = seed_post(device)

    toc1 = time.time()
    print(f'*** Model setup time: {(toc1 - tic1):.2f}s')

    counter = 0
    with torch.no_grad():
        with precision_scope(device.type):
            with model.ema_scope():

                for iteration in range(ITERATIONS):
                    for prompt, negative in get_prompts():
                        print(f'*** Iteration {iteration + 1}: {prompt}')
                        tic2 = time.time()
                        images = generate_samples(model, sampler, prompt, negative, start_code)
                        for image in images:
                            name = save_image(image)
                            save_history(name, prompt, negative)
                            print(f'*** Saved image: {name}')
                        toc2 = time.time()

                        print(f'*** Synthesis time: {(toc2 - tic2):.2f}s')
                        counter += len(images)

    print(f'*** Total time: {(toc2 - tic1):.2f}s')
    print(f'*** Saved {counter} image(s) to {FOLDER} folder.')

if __name__ == "__main__":
    except KeyboardInterrupt:
        print('*** User abort, goodbye.')
    except FileNotFoundError as e:
        print(f'*** {e}')

Aside: I found Ryan O’Connor’s article “How to Run Stable Diffusion Locally to Generate Images” on AssemblyAI to be simple enough for me to understand the Stable Diffusion options and prompt engineering. Also, ArtHub.ai has loads of sample prompts and the images generated by users.